# Stable Diffusion optimization

An anime-style text-to-image generation model based on SDXL, suitable for beginners, capable of generating high-quality anime characters and stylized images.

An openpose controlnet model trained for flux-dev, used for controlling image generation through human poses

Based on Stable Diffusion model with LoRA technology, specializing in generating anime-style images.

A LoRA fine-tuned model based on Stable Diffusion, specialized in generating retro collage-style images

A LoRA fine-tuned model based on Stable Diffusion, specifically designed for generating illustrations in the Victorian-era drawing style

A LoRA fine-tuned model based on anime style, specifically designed for generating high-quality anime-style images

A fine-tuned version based on the Stable Diffusion architecture, utilizing Low-Rank Adaptation (LoRA) technology, focusing on high-quality face image generation.

Anything V7 is a text-to-image generation model based on Stable Diffusion technology, specializing in anime, animation, and cartoon-style image generation.

CyberRealistic is a model suitable for Stable Diffusion Webui Automatic1111, which can be used to generate images with photorealistic and cyber-realistic styles.

A Stable Diffusion model tailored for generating cartoon and anime-style creative concept images

Kencanmix v2.0 Beta is a text-to-image generation model based on Stable Diffusion technology, capable of producing high-quality images from textual descriptions.

CyberRealistic_V3 is a diffusion model based on Stable Diffusion that can convert text descriptions into realistic images.

An image generation model fine-tuned based on Stable Diffusion 1.5, focusing on generating semi-realistic style images.

Stable Diffusion model fine-tuned with Dreambooth technology for generating high-quality personalized images

An image generation model fine-tuned based on Stable Diffusion 1.5, focusing on realistic and artistic style image generation.

Anything V5 is a Stable Diffusion model developed by the original creator of Anything V3, focusing on text-to-image generation and supporting high-quality artistic creation.

A text-to-image generation model based on Stable Diffusion technology, specializing in generating images related to the style of 'The Eminence in Shadow'.

The Openjourney fine-tuned model is trained based on Stable Diffusion 1.5, using over 60,000 images for training to optimize image generation quality.

This is a Dreambooth model based on Stable Diffusion, trained by n0rf using buildspace's DreamBooth notebook, primarily for generating personalized AI avatars.

The JRAD Dreambooth model is a text-to-image generation model trained by ChrtBlnch using buildspace's DreamBooth notebook, ideal for creating personalized AI avatars.

A text-to-image generation model based on Stable Diffusion, specializing in artistic style creation

This is a Stable Diffusion model trained using Dreambooth technology for text-to-image generation tasks.

An anime-style image generation model trained on the Universal Model V3.0, specifically designed to produce high-quality images of Banano (the cryptocurrency mascot).

A text-to-image generation model based on Stable Diffusion, specializing in fashion and 3D style image generation.

This is a text-to-image generation model based on Stable Diffusion technology, trained using the Dreambooth method.

A text-to-image model fine-tuned on Stable Diffusion using Dreambooth, learning the concept of noggles_fastdb_4800

Openjourney is an open-source Stable Diffusion model fine-tuned by PromptHero based on Midjourney images, specializing in generating high-quality artistic images.

A Stable Diffusion model trained with Dreambooth, specifically designed for generating 'Marco'-style female images.

A Stable Diffusion model fine-tuned on scenes from the TV series 'Arcane', capable of triggering style effects with specific tags